Abstract: Diffusion models, such as Stable Diffusion, have shown incredible performance on text-to-image generation. Since text-to-image generation often requires models to generate visual concepts with fine-grained details and attributes specified in text prompts, can we leverage the powerful representations learned by pre-trained diffusion models for discriminative tasks such as image-text matching? To answer this question, we propose a novel approach, Discriminative Stable Diffusion (Discffusion), which turns pre-trained text-to-image diffusion models into few-shot discriminative learners. Our approach uses the cross-attention score of a Stable Diffusion model to capture the mutual influence between visual and textual information and fine-tune the model via a new attention-based prompt learning to perform image-text matching. By comparing Discffusion with state-of-the-art methods on several benchmark datasets, we demonstrate the potential of using pre-trained diffusion models for discriminative tasks with superior results on few-shot image-text matching.
Submission Length: Regular submission (no more than 12 pages of main content)
Changes Since Last Submission: We have made the following updates to the paper:
- Created a new Figure 2 and updated the figure captions.
- Rewritten the method section, including the algorithm block.
- Provided additional intuitions and explanations in the method section.
- Reorganized the experiments section and added additional experiments.
- Reorganized all tables, figures, and algorithms, enhancing the overall presentation of the paper.
- Included additional experimental settings.
- Revised the conclusion section.
Overall, we have also thoroughly proofread the entire paper, improved the writing, provided additional clarifications, and included more details. All modifications are highlighted in blue.
Code: https://github.com/eric-ai-lab/Discffusion
Supplementary Material: zip
Assigned Action Editor: ~Yingzhen_Li1
Submission Number: 2556
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